Sebastian Golebiewskiandbstankix
6761a29af5
* update-notebooks * Update docs/nbdoc/nbdoc.py Co-authored-by: bstankix <bartoszx.stankiewicz@intel.com> * Update docs/nbdoc/nbdoc.py Co-authored-by: bstankix <bartoszx.stankiewicz@intel.com> --------- Co-authored-by: bstankix <bartoszx.stankiewicz@intel.com>
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1006 lines
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Text-to-Image Generation with Stable Diffusion v2 and OpenVINO™
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===============================================================
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Stable Diffusion v2 is the next generation of Stable Diffusion model a
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Text-to-Image latent diffusion model created by the researchers and
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engineers from `Stability AI <https://stability.ai/>`__ and
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`LAION <https://laion.ai/>`__.
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General diffusion models are machine learning systems that are trained
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to denoise random gaussian noise step by step, to get to a sample of
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interest, such as an image. Diffusion models have shown to achieve
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state-of-the-art results for generating image data. But one downside of
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diffusion models is that the reverse denoising process is slow. In
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addition, these models consume a lot of memory because they operate in
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pixel space, which becomes unreasonably expensive when generating
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high-resolution images. Therefore, it is challenging to train these
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models and also use them for inference. OpenVINO brings capabilities to
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run model inference on Intel hardware and opens the door to the
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fantastic world of diffusion models for everyone!
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In previous notebooks, we already discussed how to run `Text-to-Image
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generation and Image-to-Image generation using Stable Diffusion
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v1 <225-stable-diffusion-text-to-image-with-output.html>`__
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and `controlling its generation process using
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ControlNet <./235-controlnet-stable-diffusion/235-controlnet-stable-diffusion.ipynb>`__.
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Now is turn of Stable Diffusion v2.
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Stable Diffusion v2: What’s new?
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--------------------------------
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The new stable diffusion model offers a bunch of new features inspired
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by the other models that have emerged since the introduction of the
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first iteration. Some of the features that can be found in the new model
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are:
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* The model comes with a new robust encoder, OpenCLIP, created by LAION and aided by Stability AI; this version v2 significantly enhances the produced photos over the V1 versions.
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* The model can now generate images in a 768x768 resolution, offering more information to be shown in the generated images.
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* The model finetuned with `v-objective <https://arxiv.org/abs/2202.00512>`__. The v-parameterization is particularly useful for numerical stability throughout the diffusion process to enable progressive distillation for models. For models that operate at higher resolution, it is also discovered that the v-parameterization avoids color shifting artifacts that are known to affect high resolution diffusion models, and in the video setting it avoids temporal color shifting that sometimes appears with epsilon-prediction used in Stable Diffusion v1.
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* The model also comes with a new diffusion model capable of running upscaling on the images generated. Upscaled images can be adjusted up to 4 times the original image. Provided as separated model, for more details please check `stable-diffusion-x4-upscaler <https://huggingface.co/stabilityai/stable-diffusion-x4-upscaler>`__
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* The model comes with a new refined depth architecture capable of preserving context from prior generation layers in an img2img setting. This structure preservation helps generate images that preserving forms and shadow of objects, but with different content.
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* The model comes with an updated inpainting module built upon the previous model. This text-guided inpainting makes switching out parts in the image easier than before.
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This notebook demonstrates how to convert and run Stable Diffusion v2
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model using OpenVINO.
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Notebook contains the following steps:
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1. Convert PyTorch models to ONNX format.
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2. Convert ONNX models to OpenVINO IR format, using Model Optimizer tool.
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3. Run Stable Diffusion v2 Text-to-Image pipeline with OpenVINO.
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**Note:** This is the full version of the Stable Diffusion text-to-image
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implementation. If you would like to get started and run the notebook
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quickly, check out
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`236-stable-diffusion-v2-text-to-image-demo.ipynb <https://github.com/openvinotoolkit/openvino_notebooks/blob/main/notebooks/236-stable-diffusion-v2/236-stable-diffusion-v2-text-to-image-demo.ipynb>`__.
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Prerequisites
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-------------
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install required packages
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.. code:: ipython3
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!pip install -q "diffusers>=0.14.0" openvino-dev "transformers >= 4.25.1" gradio
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.. parsed-literal::
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[notice] A new release of pip available: 22.3.1 -> 23.0.1
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[notice] To update, run: pip install --upgrade pip
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Stable Diffusion v2 for Text-to-Image Generation
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------------------------------------------------
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To start, let’s look on Text-to-Image process for Stable Diffusion v2.
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We will use
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`stabilitai/stable-diffusion-2-1 <https://huggingface.co/stabilityai/stable-diffusion-2-1>`__
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model for these purposes. The main difference from Stable Diffusion v2
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and Stable Diffusion v2.1 is usage of more data, more training, and less
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restrictive filtering of the dataset, that gives promising results for
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selecting wide range of input text prompts. More details about model can
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be found in `Stability AI blog
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post <https://stability.ai/blog/stablediffusion2-1-release7-dec-2022>`__
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and original model
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`repository <https://github.com/Stability-AI/stablediffusion>`__.
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Stable Diffusion in Diffusers library
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~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~
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To work with Stable Diffusion v2, we will use Hugging Face
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`Diffusers <https://github.com/huggingface/diffusers>`__ library. To
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experiment with Stable Diffusion models, Diffusers exposes the
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`StableDiffusionPipeline <https://huggingface.co/docs/diffusers/using-diffusers/conditional_image_generation>`__
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similar to the `other Diffusers
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pipelines <https://huggingface.co/docs/diffusers/api/pipelines/overview>`__.
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The code below demonstrates how to create ``StableDiffusionPipeline``
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using ``stable-diffusion-2-1``:
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.. code:: ipython3
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from diffusers import StableDiffusionPipeline
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pipe = StableDiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2-1-base").to("cpu")
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# for reducing memory consumption get all components from pipeline independently
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text_encoder = pipe.text_encoder
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text_encoder.eval()
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unet = pipe.unet
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unet.eval()
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vae = pipe.vae
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vae.eval()
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conf = pipe.scheduler.config
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del pipe
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.. parsed-literal::
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Fetching 13 files: 0%| | 0/13 [00:00<?, ?it/s]
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.. parsed-literal::
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/home/ea/work/transformers/src/transformers/models/clip/feature_extraction_clip.py:28: FutureWarning: The class CLIPFeatureExtractor is deprecated and will be removed in version 5 of Transformers. Please use CLIPImageProcessor instead.
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warnings.warn(
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Convert models to OpenVINO Intermediate representation (IR) format
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~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~~
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OpenVINO supports PyTorch through export to the ONNX format. We will use
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the ``torch.onnx.export`` function to obtain the ONNX model, we can
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learn more in the `PyTorch
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documentation <https://pytorch.org/docs/stable/onnx.html>`__. We need to
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provide a model object, input data for model tracing, and a path for
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saving the model. Optionally, we can provide a target ONNX opset for
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conversion and other parameters specified in the documentation (for
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example, input and output names or dynamic shapes).
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While ONNX models are directly supported by OpenVINO™ runtime, it can be
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useful to convert them to IR format to take the advantage of advanced
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OpenVINO optimization tools and features. We will use OpenVINO `Model
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Optimizer <https://docs.openvino.ai/2023.0/openvino_docs_MO_DG_Deep_Learning_Model_Optimizer_DevGuide.html>`__
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to convert a model to IR format.
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The pipeline consists of three important parts:
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* Text Encoder to create condition to generate an image from a text prompt.
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* U-Net for step-by-step denoising latent image representation.
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* Autoencoder (VAE) for decoding latent space to image.
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Let us convert each part:
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Text Encoder
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~~~~~~~~~~~~
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The text-encoder is responsible for transforming the input prompt, for
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example, “a photo of an astronaut riding a horse” into an embedding
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space that can be understood by the U-Net. It is usually a simple
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transformer-based encoder that maps a sequence of input tokens to a
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sequence of latent text embeddings.
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The input of the text encoder is tensor ``input_ids``, which contains
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indexes of tokens from text processed by the tokenizer and padded to the
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maximum length accepted by the model. Model outputs are two tensors:
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``last_hidden_state`` - hidden state from the last MultiHeadAttention
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layer in the model and ``pooler_out`` - pooled output for whole model
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hidden states. We will use ``opset_version=14`` because the model
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contains the ``triu`` operation, supported in ONNX only starting from
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this opset.
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.. code:: ipython3
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from pathlib import Path
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sd2_1_model_dir = Path("sd2.1")
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sd2_1_model_dir.mkdir(exist_ok=True)
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.. code:: ipython3
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import gc
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import torch
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TEXT_ENCODER_ONNX_PATH = sd2_1_model_dir / 'text_encoder.onnx'
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TEXT_ENCODER_OV_PATH = TEXT_ENCODER_ONNX_PATH.with_suffix('.xml')
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def convert_encoder_onnx(text_encoder: torch.nn.Module, onnx_path:Path):
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"""
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Convert Text Encoder model to ONNX.
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Function accepts pipeline, prepares example inputs for ONNX conversion via torch.export,
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Parameters:
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text_encoder (torch.nn.Module): text encoder PyTorch model
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onnx_path (Path): File for storing onnx model
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Returns:
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None
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"""
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if not onnx_path.exists():
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input_ids = torch.ones((1, 77), dtype=torch.long)
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# switch model to inference mode
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text_encoder.eval()
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# disable gradients calculation for reducing memory consumption
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with torch.no_grad():
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# export model to ONNX format
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torch.onnx._export(
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text_encoder, # model instance
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input_ids, # inputs for model tracing
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onnx_path, # output file for saving result
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input_names=['tokens'], # model input name for onnx representation
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output_names=['last_hidden_state', 'pooler_out'], # model output names for onnx representation
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opset_version=14, # onnx opset version for export,
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onnx_shape_inference=False
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)
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print('Text Encoder successfully converted to ONNX')
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if not TEXT_ENCODER_OV_PATH.exists():
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convert_encoder_onnx(text_encoder, TEXT_ENCODER_ONNX_PATH)
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!mo --input_model $TEXT_ENCODER_ONNX_PATH --output_dir $sd2_1_model_dir
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print('Text Encoder successfully converted to IR')
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else:
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print(f"Text encoder will be loaded from {TEXT_ENCODER_OV_PATH}")
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del text_encoder
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gc.collect();
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.. parsed-literal::
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/tmp/ipykernel_383583/1233802758.py:26: FutureWarning: 'torch.onnx._export' is deprecated in version 1.12.0 and will be removed in version 1.14. Please use `torch.onnx.export` instead.
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torch.onnx._export(
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/home/ea/work/transformers/src/transformers/models/clip/modeling_clip.py:759: TracerWarning: torch.tensor results are registered as constants in the trace. You can safely ignore this warning if you use this function to create tensors out of constant variables that would be the same every time you call this function. In any other case, this might cause the trace to be incorrect.
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mask.fill_(torch.tensor(torch.finfo(dtype).min))
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/home/ea/work/transformers/src/transformers/models/clip/modeling_clip.py:284: TracerWarning: Converting a tensor to a Python boolean might cause the trace to be incorrect. We can't record the data flow of Python values, so this value will be treated as a constant in the future. This means that the trace might not generalize to other inputs!
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if attn_weights.size() != (bsz * self.num_heads, tgt_len, src_len):
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/home/ea/work/transformers/src/transformers/models/clip/modeling_clip.py:292: TracerWarning: Converting a tensor to a Python boolean might cause the trace to be incorrect. We can't record the data flow of Python values, so this value will be treated as a constant in the future. This means that the trace might not generalize to other inputs!
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if causal_attention_mask.size() != (bsz, 1, tgt_len, src_len):
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/home/ea/work/transformers/src/transformers/models/clip/modeling_clip.py:324: TracerWarning: Converting a tensor to a Python boolean might cause the trace to be incorrect. We can't record the data flow of Python values, so this value will be treated as a constant in the future. This means that the trace might not generalize to other inputs!
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if attn_output.size() != (bsz * self.num_heads, tgt_len, self.head_dim):
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/home/ea/work/notebooks_env/lib/python3.8/site-packages/torch/onnx/symbolic_helper.py:710: UserWarning: Type cannot be inferred, which might cause exported graph to produce incorrect results.
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warnings.warn(
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/home/ea/work/notebooks_env/lib/python3.8/site-packages/torch/onnx/symbolic_opset9.py:5408: UserWarning: Exporting aten::index operator of advanced indexing in opset 14 is achieved by combination of multiple ONNX operators, including Reshape, Transpose, Concat, and Gather. If indices include negative values, the exported graph will produce incorrect results.
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warnings.warn(
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.. parsed-literal::
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Text Encoder successfully converted to ONNX
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Warning: One or more of the values of the Constant can't fit in the float16 data type. Those values were casted to the nearest limit value, the model can produce incorrect results.
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[ INFO ] The model was converted to IR v11, the latest model format that corresponds to the source DL framework input/output format. While IR v11 is backwards compatible with OpenVINO Inference Engine API v1.0, please use API v2.0 (as of 2022.1) to take advantage of the latest improvements in IR v11.
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Find more information about API v2.0 and IR v11 at https://docs.openvino.ai/2023.0/openvino_2_0_transition_guide.html
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[ SUCCESS ] Generated IR version 11 model.
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[ SUCCESS ] XML file: /home/ea/work/openvino_notebooks/notebooks/236-stable-diffusion-v2/sd2.1/text_encoder.xml
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[ SUCCESS ] BIN file: /home/ea/work/openvino_notebooks/notebooks/236-stable-diffusion-v2/sd2.1/text_encoder.bin
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Text Encoder successfully converted to IR
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U-Net
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~~~~~
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U-Net model gradually denoises latent image representation guided by
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text encoder hidden state.
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U-Net model has three inputs:
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* ``sample`` - latent image sample from previous step. Generation process has not been started yet, so you will use random noise.
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* ``timestep`` - current scheduler step.
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* ``encoder_hidden_state`` - hidden state of text encoder.
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Model predicts the ``sample`` state for the next step.
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Generally, U-Net model conversion process remain the same like in Stable
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Diffusion v1, expect small changes in input sample size. Our model was
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pretrained to generate images with resolution 768x768, initial latent
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sample size for this case is 96x96. Besides that, for different use
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cases like inpainting and depth to image generation model also can
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accept additional image information: depth map or mask as channel-wise
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concatenation with initial latent sample. For convering U-Net model for
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such use cases required to modify number of input channels.
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.. code:: ipython3
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import numpy as np
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UNET_ONNX_PATH = sd2_1_model_dir / 'unet/unet.onnx'
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UNET_OV_PATH = UNET_ONNX_PATH.parents[1] / 'unet.xml'
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def convert_unet_onnx(unet:torch.nn.Module, onnx_path:Path, num_channels:int = 4, width:int = 64, height:int = 64):
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"""
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Convert Unet model to ONNX, then IR format.
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Function accepts pipeline, prepares example inputs for ONNX conversion via torch.export,
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Parameters:
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unet (torch.nn.Module): UNet PyTorch model
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onnx_path (Path): File for storing onnx model
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num_channels (int, optional, 4): number of input channels
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width (int, optional, 64): input width
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height (int, optional, 64): input height
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Returns:
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None
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"""
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if not onnx_path.exists():
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# prepare inputs
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encoder_hidden_state = torch.ones((2, 77, 1024))
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latents_shape = (2, num_channels, width, height)
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latents = torch.randn(latents_shape)
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t = torch.from_numpy(np.array(1, dtype=np.float32))
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# model size > 2Gb, it will be represented as onnx with external data files, we will store it in separated directory for avoid a lot of files in current directory
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onnx_path.parent.mkdir(exist_ok=True, parents=True)
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unet.eval()
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with torch.no_grad():
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torch.onnx._export(
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unet,
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(latents, t, encoder_hidden_state), str(onnx_path),
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input_names=['latent_model_input', 't', 'encoder_hidden_states'],
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output_names=['out_sample'],
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onnx_shape_inference=False
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)
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print('U-Net successfully converted to ONNX')
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if not UNET_OV_PATH.exists():
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convert_unet_onnx(unet, UNET_ONNX_PATH, width=96, height=96)
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del unet
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gc.collect()
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!mo --input_model $UNET_ONNX_PATH --output_dir $sd2_1_model_dir
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print('U-Net successfully converted to IR')
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else:
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del unet
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print(f"U-Net will be loaded from {UNET_OV_PATH}")
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gc.collect();
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.. parsed-literal::
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/tmp/ipykernel_383583/4211352295.py:32: FutureWarning: 'torch.onnx._export' is deprecated in version 1.12.0 and will be removed in version 1.14. Please use `torch.onnx.export` instead.
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torch.onnx._export(
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/home/ea/work/diffusers/src/diffusers/models/unet_2d_condition.py:526: TracerWarning: Converting a tensor to a Python boolean might cause the trace to be incorrect. We can't record the data flow of Python values, so this value will be treated as a constant in the future. This means that the trace might not generalize to other inputs!
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if any(s % default_overall_up_factor != 0 for s in sample.shape[-2:]):
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/home/ea/work/diffusers/src/diffusers/models/resnet.py:185: TracerWarning: Converting a tensor to a Python boolean might cause the trace to be incorrect. We can't record the data flow of Python values, so this value will be treated as a constant in the future. This means that the trace might not generalize to other inputs!
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assert hidden_states.shape[1] == self.channels
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/home/ea/work/diffusers/src/diffusers/models/resnet.py:190: TracerWarning: Converting a tensor to a Python boolean might cause the trace to be incorrect. We can't record the data flow of Python values, so this value will be treated as a constant in the future. This means that the trace might not generalize to other inputs!
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assert hidden_states.shape[1] == self.channels
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/home/ea/work/diffusers/src/diffusers/models/resnet.py:112: TracerWarning: Converting a tensor to a Python boolean might cause the trace to be incorrect. We can't record the data flow of Python values, so this value will be treated as a constant in the future. This means that the trace might not generalize to other inputs!
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assert hidden_states.shape[1] == self.channels
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/home/ea/work/diffusers/src/diffusers/models/resnet.py:125: TracerWarning: Converting a tensor to a Python boolean might cause the trace to be incorrect. We can't record the data flow of Python values, so this value will be treated as a constant in the future. This means that the trace might not generalize to other inputs!
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if hidden_states.shape[0] >= 64:
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/home/ea/work/diffusers/src/diffusers/models/unet_2d_condition.py:651: TracerWarning: Converting a tensor to a Python boolean might cause the trace to be incorrect. We can't record the data flow of Python values, so this value will be treated as a constant in the future. This means that the trace might not generalize to other inputs!
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if not return_dict:
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.. parsed-literal::
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U-Net successfully converted to ONNX
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[ INFO ] The model was converted to IR v11, the latest model format that corresponds to the source DL framework input/output format. While IR v11 is backwards compatible with OpenVINO Inference Engine API v1.0, please use API v2.0 (as of 2022.1) to take advantage of the latest improvements in IR v11.
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Find more information about API v2.0 and IR v11 at https://docs.openvino.ai/2023.0/openvino_2_0_transition_guide.html
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[ SUCCESS ] Generated IR version 11 model.
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[ SUCCESS ] XML file: /home/ea/work/openvino_notebooks/notebooks/236-stable-diffusion-v2/sd2.1/unet.xml
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[ SUCCESS ] BIN file: /home/ea/work/openvino_notebooks/notebooks/236-stable-diffusion-v2/sd2.1/unet.bin
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U-Net successfully converted to IR
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VAE
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~~~
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The VAE model has two parts, an encoder and a decoder. The encoder is
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used to convert the image into a low dimensional latent representation,
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which will serve as the input to the U-Net model. The decoder,
|
||
conversely, transforms the latent representation back into an image.
|
||
|
||
During latent diffusion training, the encoder is used to get the latent
|
||
representations (latents) of the images for the forward diffusion
|
||
process, which applies more and more noise at each step. During
|
||
inference, the denoised latents generated by the reverse diffusion
|
||
process are converted back into images using the VAE decoder. When you
|
||
run inference for Text-to-Image, there is no initial image as a starting
|
||
point. You can skip this step and directly generate initial random
|
||
noise.
|
||
|
||
When running Text-to-Image pipeline, we will see that we **only need the
|
||
VAE decoder**, but preserve VAE encoder conversion, it will be useful in
|
||
next chapter of our tutorial.
|
||
|
||
Note: This process will take a few minutes and use significant amount of
|
||
RAM (recommended at least 32GB).
|
||
|
||
.. code:: ipython3
|
||
|
||
VAE_ENCODER_ONNX_PATH = sd2_1_model_dir / 'vae_encoder.onnx'
|
||
VAE_ENCODER_OV_PATH = VAE_ENCODER_ONNX_PATH.with_suffix('.xml')
|
||
|
||
|
||
def convert_vae_encoder_onnx(vae: torch.nn.Module, onnx_path: Path, width:int = 512, height:int = 512):
|
||
"""
|
||
Convert VAE model to ONNX, then IR format.
|
||
Function accepts pipeline, creates wrapper class for export only necessary for inference part,
|
||
prepares example inputs for ONNX conversion via torch.export,
|
||
Parameters:
|
||
vae (torch.nn.Module): VAE PyTorch model
|
||
onnx_path (Path): File for storing onnx model
|
||
width (int, optional, 512): input width
|
||
height (int, optional, 512): input height
|
||
Returns:
|
||
None
|
||
"""
|
||
class VAEEncoderWrapper(torch.nn.Module):
|
||
def __init__(self, vae):
|
||
super().__init__()
|
||
self.vae = vae
|
||
|
||
def forward(self, image):
|
||
h = self.vae.encoder(image)
|
||
moments = self.vae.quant_conv(h)
|
||
return moments
|
||
|
||
if not onnx_path.exists():
|
||
vae_encoder = VAEEncoderWrapper(vae)
|
||
vae_encoder.eval()
|
||
image = torch.zeros((1, 3, width, height))
|
||
with torch.no_grad():
|
||
torch.onnx.export(vae_encoder, image, onnx_path, input_names=[
|
||
'init_image'], output_names=['image_latent'])
|
||
print('VAE encoder successfully converted to ONNX')
|
||
|
||
|
||
if not VAE_ENCODER_OV_PATH.exists():
|
||
convert_vae_encoder_onnx(vae, VAE_ENCODER_ONNX_PATH, 768, 768)
|
||
!mo --input_model $VAE_ENCODER_ONNX_PATH --output_dir $sd2_1_model_dir
|
||
print('VAE encoder successfully converted to IR')
|
||
else:
|
||
print(f"VAE encoder will be loaded from {VAE_ENCODER_OV_PATH}")
|
||
|
||
VAE_DECODER_ONNX_PATH = sd2_1_model_dir / 'vae_decoder.onnx'
|
||
VAE_DECODER_OV_PATH = VAE_DECODER_ONNX_PATH.with_suffix('.xml')
|
||
|
||
|
||
def convert_vae_decoder_onnx(vae: torch.nn.Module, onnx_path: Path, width:int = 64, height:int = 64):
|
||
"""
|
||
Convert VAE model to ONNX, then IR format.
|
||
Function accepts pipeline, creates wrapper class for export only necessary for inference part,
|
||
prepares example inputs for ONNX conversion via torch.export,
|
||
Parameters:
|
||
vae:
|
||
onnx_path (Path): File for storing onnx model
|
||
width (int, optional, 64): input width
|
||
height (int, optional, 64): input height
|
||
Returns:
|
||
None
|
||
"""
|
||
class VAEDecoderWrapper(torch.nn.Module):
|
||
def __init__(self, vae):
|
||
super().__init__()
|
||
self.vae = vae
|
||
|
||
def forward(self, latents):
|
||
latents = 1 / 0.18215 * latents
|
||
return self.vae.decode(latents)
|
||
|
||
if not onnx_path.exists():
|
||
vae_decoder = VAEDecoderWrapper(vae)
|
||
latents = torch.zeros((1, 4, width, height))
|
||
|
||
vae_decoder.eval()
|
||
with torch.no_grad():
|
||
torch.onnx.export(vae_decoder, latents, onnx_path, input_names=[
|
||
'latents'], output_names=['sample'])
|
||
print('VAE decoder successfully converted to ONNX')
|
||
|
||
|
||
if not VAE_DECODER_OV_PATH.exists():
|
||
convert_vae_decoder_onnx(vae, VAE_DECODER_ONNX_PATH, 96, 96)
|
||
!mo --input_model $VAE_DECODER_ONNX_PATH --output_dir $sd2_1_model_dir
|
||
print('VAE decoder successfully converted to IR')
|
||
else:
|
||
print(f"VAE decoder will be loaded from {VAE_DECODER_OV_PATH}")
|
||
|
||
del vae
|
||
gc.collect();
|
||
|
||
|
||
.. parsed-literal::
|
||
|
||
/home/ea/work/notebooks_env/lib/python3.8/site-packages/torch/onnx/_internal/jit_utils.py:258: UserWarning: Constant folding - Only steps=1 can be constant folded for opset >= 10 onnx::Slice op. Constant folding not applied. (Triggered internally at ../torch/csrc/jit/passes/onnx/constant_fold.cpp:179.)
|
||
_C._jit_pass_onnx_node_shape_type_inference(node, params_dict, opset_version)
|
||
/home/ea/work/notebooks_env/lib/python3.8/site-packages/torch/onnx/utils.py:687: UserWarning: Constant folding - Only steps=1 can be constant folded for opset >= 10 onnx::Slice op. Constant folding not applied. (Triggered internally at ../torch/csrc/jit/passes/onnx/constant_fold.cpp:179.)
|
||
_C._jit_pass_onnx_graph_shape_type_inference(
|
||
/home/ea/work/notebooks_env/lib/python3.8/site-packages/torch/onnx/utils.py:1178: UserWarning: Constant folding - Only steps=1 can be constant folded for opset >= 10 onnx::Slice op. Constant folding not applied. (Triggered internally at ../torch/csrc/jit/passes/onnx/constant_fold.cpp:179.)
|
||
_C._jit_pass_onnx_graph_shape_type_inference(
|
||
|
||
|
||
.. parsed-literal::
|
||
|
||
VAE encoder successfully converted to ONNX
|
||
[ INFO ] The model was converted to IR v11, the latest model format that corresponds to the source DL framework input/output format. While IR v11 is backwards compatible with OpenVINO Inference Engine API v1.0, please use API v2.0 (as of 2022.1) to take advantage of the latest improvements in IR v11.
|
||
Find more information about API v2.0 and IR v11 at https://docs.openvino.ai/2023.0/openvino_2_0_transition_guide.html
|
||
[ SUCCESS ] Generated IR version 11 model.
|
||
[ SUCCESS ] XML file: /home/ea/work/openvino_notebooks/notebooks/236-stable-diffusion-v2/sd2.1/vae_encoder.xml
|
||
[ SUCCESS ] BIN file: /home/ea/work/openvino_notebooks/notebooks/236-stable-diffusion-v2/sd2.1/vae_encoder.bin
|
||
VAE encoder successfully converted to IR
|
||
|
||
|
||
.. parsed-literal::
|
||
|
||
/home/ea/work/notebooks_env/lib/python3.8/site-packages/torch/onnx/_internal/jit_utils.py:258: UserWarning: The shape inference of prim::Constant type is missing, so it may result in wrong shape inference for the exported graph. Please consider adding it in symbolic function. (Triggered internally at ../torch/csrc/jit/passes/onnx/shape_type_inference.cpp:1884.)
|
||
_C._jit_pass_onnx_node_shape_type_inference(node, params_dict, opset_version)
|
||
/home/ea/work/notebooks_env/lib/python3.8/site-packages/torch/onnx/utils.py:687: UserWarning: The shape inference of prim::Constant type is missing, so it may result in wrong shape inference for the exported graph. Please consider adding it in symbolic function. (Triggered internally at ../torch/csrc/jit/passes/onnx/shape_type_inference.cpp:1884.)
|
||
_C._jit_pass_onnx_graph_shape_type_inference(
|
||
/home/ea/work/notebooks_env/lib/python3.8/site-packages/torch/onnx/utils.py:1178: UserWarning: The shape inference of prim::Constant type is missing, so it may result in wrong shape inference for the exported graph. Please consider adding it in symbolic function. (Triggered internally at ../torch/csrc/jit/passes/onnx/shape_type_inference.cpp:1884.)
|
||
_C._jit_pass_onnx_graph_shape_type_inference(
|
||
|
||
|
||
.. parsed-literal::
|
||
|
||
VAE decoder successfully converted to ONNX
|
||
[ INFO ] The model was converted to IR v11, the latest model format that corresponds to the source DL framework input/output format. While IR v11 is backwards compatible with OpenVINO Inference Engine API v1.0, please use API v2.0 (as of 2022.1) to take advantage of the latest improvements in IR v11.
|
||
Find more information about API v2.0 and IR v11 at https://docs.openvino.ai/2023.0/openvino_2_0_transition_guide.html
|
||
[ SUCCESS ] Generated IR version 11 model.
|
||
[ SUCCESS ] XML file: /home/ea/work/openvino_notebooks/notebooks/236-stable-diffusion-v2/sd2.1/vae_decoder.xml
|
||
[ SUCCESS ] BIN file: /home/ea/work/openvino_notebooks/notebooks/236-stable-diffusion-v2/sd2.1/vae_decoder.bin
|
||
VAE decoder successfully converted to IR
|
||
|
||
|
||
Prepare Inference Pipeline
|
||
~~~~~~~~~~~~~~~~~~~~~~~~~~
|
||
|
||
Putting it all together, let us now take a closer look at how the model
|
||
works in inference by illustrating the logical flow.
|
||
|
||
.. figure:: https://github.com/openvinotoolkit/openvino_notebooks/assets/22090501/ec454103-0d28-48e3-a18e-b55da3fab381
|
||
:alt: text2img-stable-diffusion v2
|
||
|
||
text2img-stable-diffusion v2
|
||
|
||
The stable diffusion model takes both a latent seed and a text prompt as
|
||
input. The latent seed is then used to generate random latent image
|
||
representations of size :math:`96 \times 96` where as the text prompt is
|
||
transformed to text embeddings of size :math:`77 \times 1024` via
|
||
OpenCLIP’s text encoder.
|
||
|
||
Next, the U-Net iteratively *denoises* the random latent image
|
||
representations while being conditioned on the text embeddings. The
|
||
output of the U-Net, being the noise residual, is used to compute a
|
||
denoised latent image representation via a scheduler algorithm. Many
|
||
different scheduler algorithms can be used for this computation, each
|
||
having its pros and cons. For Stable Diffusion, it is recommended to use
|
||
one of:
|
||
|
||
- `PNDM
|
||
scheduler <https://github.com/huggingface/diffusers/blob/main/src/diffusers/schedulers/scheduling_pndm.py>`__
|
||
- `DDIM
|
||
scheduler <https://github.com/huggingface/diffusers/blob/main/src/diffusers/schedulers/scheduling_ddim.py>`__
|
||
- `K-LMS
|
||
scheduler <https://github.com/huggingface/diffusers/blob/main/src/diffusers/schedulers/scheduling_lms_discrete.py>`__
|
||
|
||
Theory on how the scheduler algorithm function works is out of scope for
|
||
this notebook, but in short, you should remember that they compute the
|
||
predicted denoised image representation from the previous noise
|
||
representation and the predicted noise residual. For more information,
|
||
it is recommended to look into `Elucidating the Design Space of
|
||
Diffusion-Based Generative Models <https://arxiv.org/abs/2206.00364>`__.
|
||
|
||
The chart above looks very similar to Stable Diffusion V1 from
|
||
`notebook <225-stable-diffusion-text-to-image-with-output.html>`__,
|
||
but there is some small difference in details:
|
||
|
||
* Changed input resolution for U-Net model.
|
||
* Changed text encoder and as the result size of its hidden state embeddings.
|
||
* Additionally, to improve image generation quality authors introduced negative prompting. Technically, positive prompt steers the diffusion toward the images associated with it, while negative prompt steers the diffusion away from it.In other words, negative prompt declares undesired concepts for generation image, e.g. if we want to have colorful and bright image, gray scale image will be result which we want to avoid, in this case grey scale can be treated as negative prompt. The positive and negative prompt are in equal footing. You can always use one with or without the other. More explanation of how it works can be found in this `article <https://stable-diffusion-art.com/how-negative-prompt-work/>`__.
|
||
|
||
.. code:: ipython3
|
||
|
||
import inspect
|
||
from typing import List, Optional, Union, Dict
|
||
|
||
import PIL
|
||
import cv2
|
||
import torch
|
||
|
||
from transformers import CLIPTokenizer
|
||
from diffusers.pipeline_utils import DiffusionPipeline
|
||
from diffusers.schedulers import DDIMScheduler, LMSDiscreteScheduler, PNDMScheduler
|
||
from openvino.runtime import Model
|
||
|
||
|
||
def scale_fit_to_window(dst_width:int, dst_height:int, image_width:int, image_height:int):
|
||
"""
|
||
Preprocessing helper function for calculating image size for resize with peserving original aspect ratio
|
||
and fitting image to specific window size
|
||
|
||
Parameters:
|
||
dst_width (int): destination window width
|
||
dst_height (int): destination window height
|
||
image_width (int): source image width
|
||
image_height (int): source image height
|
||
Returns:
|
||
result_width (int): calculated width for resize
|
||
result_height (int): calculated height for resize
|
||
"""
|
||
im_scale = min(dst_height / image_height, dst_width / image_width)
|
||
return int(im_scale * image_width), int(im_scale * image_height)
|
||
|
||
|
||
def preprocess(image: PIL.Image.Image):
|
||
"""
|
||
Image preprocessing function. Takes image in PIL.Image format, resizes it to keep aspect ration and fits to model input window 512x512,
|
||
then converts it to np.ndarray and adds padding with zeros on right or bottom side of image (depends from aspect ratio), after that
|
||
converts data to float32 data type and change range of values from [0, 255] to [-1, 1], finally, converts data layout from planar NHWC to NCHW.
|
||
The function returns preprocessed input tensor and padding size, which can be used in postprocessing.
|
||
|
||
Parameters:
|
||
image (PIL.Image.Image): input image
|
||
Returns:
|
||
image (np.ndarray): preprocessed image tensor
|
||
meta (Dict): dictionary with preprocessing metadata info
|
||
"""
|
||
src_width, src_height = image.size
|
||
dst_width, dst_height = scale_fit_to_window(
|
||
512, 512, src_width, src_height)
|
||
image = np.array(image.resize((dst_width, dst_height),
|
||
resample=PIL.Image.Resampling.LANCZOS))[None, :]
|
||
pad_width = 512 - dst_width
|
||
pad_height = 512 - dst_height
|
||
pad = ((0, 0), (0, pad_height), (0, pad_width), (0, 0))
|
||
image = np.pad(image, pad, mode="constant")
|
||
image = image.astype(np.float32) / 255.0
|
||
image = 2.0 * image - 1.0
|
||
image = image.transpose(0, 3, 1, 2)
|
||
return image, {"padding": pad, "src_width": src_width, "src_height": src_height}
|
||
|
||
|
||
class OVStableDiffusionPipeline(DiffusionPipeline):
|
||
def __init__(
|
||
self,
|
||
vae_decoder: Model,
|
||
text_encoder: Model,
|
||
tokenizer: CLIPTokenizer,
|
||
unet: Model,
|
||
scheduler: Union[DDIMScheduler, PNDMScheduler, LMSDiscreteScheduler],
|
||
vae_encoder: Model = None,
|
||
):
|
||
"""
|
||
Pipeline for text-to-image generation using Stable Diffusion.
|
||
Parameters:
|
||
vae_decoder (Model):
|
||
Variational Auto-Encoder (VAE) Model to decode images to and from latent representations.
|
||
text_encoder (Model):
|
||
Frozen text-encoder. Stable Diffusion uses the text portion of
|
||
[CLIP](https://huggingface.co/docs/transformers/model_doc/clip#transformers.CLIPTextModel), specifically
|
||
the clip-vit-large-patch14(https://huggingface.co/openai/clip-vit-large-patch14) variant.
|
||
tokenizer (CLIPTokenizer):
|
||
Tokenizer of class CLIPTokenizer(https://huggingface.co/docs/transformers/v4.21.0/en/model_doc/clip#transformers.CLIPTokenizer).
|
||
unet (Model): Conditional U-Net architecture to denoise the encoded image latents.
|
||
vae_encoder (Model):
|
||
Variational Auto-Encoder (VAE) Model to encode images to latent representation.
|
||
scheduler (SchedulerMixin):
|
||
A scheduler to be used in combination with unet to denoise the encoded image latents. Can be one of
|
||
DDIMScheduler, LMSDiscreteScheduler, or PNDMScheduler.
|
||
"""
|
||
super().__init__()
|
||
self.scheduler = scheduler
|
||
self.vae_decoder = vae_decoder
|
||
self.vae_encoder = vae_encoder
|
||
self.text_encoder = text_encoder
|
||
self.unet = unet
|
||
self._text_encoder_output = text_encoder.output(0)
|
||
self._unet_output = unet.output(0)
|
||
self._vae_d_output = vae_decoder.output(0)
|
||
self._vae_e_output = vae_encoder.output(0) if vae_encoder is not None else None
|
||
self.height = self.unet.input(0).shape[2] * 8
|
||
self.width = self.unet.input(0).shape[3] * 8
|
||
self.tokenizer = tokenizer
|
||
|
||
def __call__(
|
||
self,
|
||
prompt: Union[str, List[str]],
|
||
image: PIL.Image.Image = None,
|
||
negative_prompt: Union[str, List[str]] = None,
|
||
num_inference_steps: Optional[int] = 50,
|
||
guidance_scale: Optional[float] = 7.5,
|
||
eta: Optional[float] = 0.0,
|
||
output_type: Optional[str] = "pil",
|
||
seed: Optional[int] = None,
|
||
strength: float = 1.0,
|
||
):
|
||
"""
|
||
Function invoked when calling the pipeline for generation.
|
||
Parameters:
|
||
prompt (str or List[str]):
|
||
The prompt or prompts to guide the image generation.
|
||
image (PIL.Image.Image, *optional*, None):
|
||
Intinal image for generation.
|
||
negative_prompt (str or List[str]):
|
||
The negative prompt or prompts to guide the image generation.
|
||
num_inference_steps (int, *optional*, defaults to 50):
|
||
The number of denoising steps. More denoising steps usually lead to a higher quality image at the
|
||
expense of slower inference.
|
||
guidance_scale (float, *optional*, defaults to 7.5):
|
||
Guidance scale as defined in Classifier-Free Diffusion Guidance(https://arxiv.org/abs/2207.12598).
|
||
guidance_scale is defined as `w` of equation 2.
|
||
Higher guidance scale encourages to generate images that are closely linked to the text prompt,
|
||
usually at the expense of lower image quality.
|
||
eta (float, *optional*, defaults to 0.0):
|
||
Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to
|
||
[DDIMScheduler], will be ignored for others.
|
||
output_type (`str`, *optional*, defaults to "pil"):
|
||
The output format of the generate image. Choose between
|
||
[PIL](https://pillow.readthedocs.io/en/stable/): PIL.Image.Image or np.array.
|
||
seed (int, *optional*, None):
|
||
Seed for random generator state initialization.
|
||
strength (int, *optional*, 1.0):
|
||
strength between initial image and generated in Image-to-Image pipeline, do not used in Text-to-Image
|
||
Returns:
|
||
Dictionary with keys:
|
||
sample - the last generated image PIL.Image.Image or np.array
|
||
"""
|
||
if seed is not None:
|
||
np.random.seed(seed)
|
||
# here `guidance_scale` is defined analog to the guidance weight `w` of equation (2)
|
||
# of the Imagen paper: https://arxiv.org/pdf/2205.11487.pdf . `guidance_scale = 1`
|
||
# corresponds to doing no classifier free guidance.
|
||
do_classifier_free_guidance = guidance_scale > 1.0
|
||
# get prompt text embeddings
|
||
text_embeddings = self._encode_prompt(prompt, do_classifier_free_guidance=do_classifier_free_guidance, negative_prompt=negative_prompt)
|
||
# set timesteps
|
||
accepts_offset = "offset" in set(inspect.signature(self.scheduler.set_timesteps).parameters.keys())
|
||
extra_set_kwargs = {}
|
||
if accepts_offset:
|
||
extra_set_kwargs["offset"] = 1
|
||
|
||
self.scheduler.set_timesteps(num_inference_steps, **extra_set_kwargs)
|
||
timesteps, num_inference_steps = self.get_timesteps(num_inference_steps, strength)
|
||
latent_timestep = timesteps[:1]
|
||
|
||
# get the initial random noise unless the user supplied it
|
||
latents, meta = self.prepare_latents(image, latent_timestep)
|
||
|
||
# prepare extra kwargs for the scheduler step, since not all schedulers have the same signature
|
||
# eta (η) is only used with the DDIMScheduler, it will be ignored for other schedulers.
|
||
# eta corresponds to η in DDIM paper: https://arxiv.org/abs/2010.02502
|
||
# and should be between [0, 1]
|
||
accepts_eta = "eta" in set(inspect.signature(self.scheduler.step).parameters.keys())
|
||
extra_step_kwargs = {}
|
||
if accepts_eta:
|
||
extra_step_kwargs["eta"] = eta
|
||
|
||
for t in self.progress_bar(timesteps):
|
||
# expand the latents if we are doing classifier free guidance
|
||
latent_model_input = np.concatenate([latents] * 2) if do_classifier_free_guidance else latents
|
||
latent_model_input = self.scheduler.scale_model_input(latent_model_input, t)
|
||
|
||
# predict the noise residual
|
||
noise_pred = self.unet([latent_model_input, np.array(t, dtype=np.float32), text_embeddings])[self._unet_output]
|
||
# perform guidance
|
||
if do_classifier_free_guidance:
|
||
noise_pred_uncond, noise_pred_text = noise_pred[0], noise_pred[1]
|
||
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)
|
||
|
||
# compute the previous noisy sample x_t -> x_t-1
|
||
latents = self.scheduler.step(torch.from_numpy(noise_pred), t, torch.from_numpy(latents), **extra_step_kwargs)["prev_sample"].numpy()
|
||
# scale and decode the image latents with vae
|
||
image = self.vae_decoder(latents)[self._vae_d_output]
|
||
|
||
image = self.postprocess_image(image, meta, output_type)
|
||
return {"sample": image}
|
||
|
||
def _encode_prompt(self, prompt:Union[str, List[str]], num_images_per_prompt:int = 1, do_classifier_free_guidance:bool = True, negative_prompt:Union[str, List[str]] = None):
|
||
"""
|
||
Encodes the prompt into text encoder hidden states.
|
||
|
||
Parameters:
|
||
prompt (str or list(str)): prompt to be encoded
|
||
num_images_per_prompt (int): number of images that should be generated per prompt
|
||
do_classifier_free_guidance (bool): whether to use classifier free guidance or not
|
||
negative_prompt (str or list(str)): negative prompt to be encoded
|
||
Returns:
|
||
text_embeddings (np.ndarray): text encoder hidden states
|
||
"""
|
||
batch_size = len(prompt) if isinstance(prompt, list) else 1
|
||
|
||
# tokenize input prompts
|
||
text_inputs = self.tokenizer(
|
||
prompt,
|
||
padding="max_length",
|
||
max_length=self.tokenizer.model_max_length,
|
||
truncation=True,
|
||
return_tensors="np",
|
||
)
|
||
text_input_ids = text_inputs.input_ids
|
||
|
||
text_embeddings = self.text_encoder(
|
||
text_input_ids)[self._text_encoder_output]
|
||
|
||
# duplicate text embeddings for each generation per prompt
|
||
if num_images_per_prompt != 1:
|
||
bs_embed, seq_len, _ = text_embeddings.shape
|
||
text_embeddings = np.tile(
|
||
text_embeddings, (1, num_images_per_prompt, 1))
|
||
text_embeddings = np.reshape(
|
||
text_embeddings, (bs_embed * num_images_per_prompt, seq_len, -1))
|
||
|
||
# get unconditional embeddings for classifier free guidance
|
||
if do_classifier_free_guidance:
|
||
uncond_tokens: List[str]
|
||
max_length = text_input_ids.shape[-1]
|
||
if negative_prompt is None:
|
||
uncond_tokens = [""] * batch_size
|
||
elif isinstance(negative_prompt, str):
|
||
uncond_tokens = [negative_prompt]
|
||
else:
|
||
uncond_tokens = negative_prompt
|
||
uncond_input = self.tokenizer(
|
||
uncond_tokens,
|
||
padding="max_length",
|
||
max_length=max_length,
|
||
truncation=True,
|
||
return_tensors="np",
|
||
)
|
||
|
||
uncond_embeddings = self.text_encoder(uncond_input.input_ids)[self._text_encoder_output]
|
||
|
||
# duplicate unconditional embeddings for each generation per prompt, using mps friendly method
|
||
seq_len = uncond_embeddings.shape[1]
|
||
uncond_embeddings = np.tile(uncond_embeddings, (1, num_images_per_prompt, 1))
|
||
uncond_embeddings = np.reshape(uncond_embeddings, (batch_size * num_images_per_prompt, seq_len, -1))
|
||
|
||
# For classifier free guidance, we need to do two forward passes.
|
||
# Here we concatenate the unconditional and text embeddings into a single batch
|
||
# to avoid doing two forward passes
|
||
text_embeddings = np.concatenate([uncond_embeddings, text_embeddings])
|
||
|
||
return text_embeddings
|
||
|
||
def prepare_latents(self, image:PIL.Image.Image = None, latent_timestep:torch.Tensor = None):
|
||
"""
|
||
Function for getting initial latents for starting generation
|
||
|
||
Parameters:
|
||
image (PIL.Image.Image, *optional*, None):
|
||
Input image for generation, if not provided randon noise will be used as starting point
|
||
latent_timestep (torch.Tensor, *optional*, None):
|
||
Predicted by scheduler initial step for image generation, required for latent image mixing with nosie
|
||
Returns:
|
||
latents (np.ndarray):
|
||
Image encoded in latent space
|
||
"""
|
||
latents_shape = (1, 4, self.height // 8, self.width // 8)
|
||
noise = np.random.randn(*latents_shape).astype(np.float32)
|
||
if image is None:
|
||
# if we use LMSDiscreteScheduler, let's make sure latents are mulitplied by sigmas
|
||
if isinstance(self.scheduler, LMSDiscreteScheduler):
|
||
noise = noise * self.scheduler.sigmas[0].numpy()
|
||
return noise, {}
|
||
input_image, meta = preprocess(image)
|
||
moments = self.vae_encoder(input_image)[self._vae_e_output]
|
||
mean, logvar = np.split(moments, 2, axis=1)
|
||
std = np.exp(logvar * 0.5)
|
||
latents = (mean + std * np.random.randn(*mean.shape)) * 0.18215
|
||
latents = self.scheduler.add_noise(torch.from_numpy(latents), torch.from_numpy(noise), latent_timestep).numpy()
|
||
return latents, meta
|
||
|
||
def postprocess_image(self, image:np.ndarray, meta:Dict, output_type:str = "pil"):
|
||
"""
|
||
Postprocessing for decoded image. Takes generated image decoded by VAE decoder, unpad it to initila image size (if required),
|
||
normalize and convert to [0, 255] pixels range. Optionally, convertes it from np.ndarray to PIL.Image format
|
||
|
||
Parameters:
|
||
image (np.ndarray):
|
||
Generated image
|
||
meta (Dict):
|
||
Metadata obtained on latents preparing step, can be empty
|
||
output_type (str, *optional*, pil):
|
||
Output format for result, can be pil or numpy
|
||
Returns:
|
||
image (List of np.ndarray or PIL.Image.Image):
|
||
Postprocessed images
|
||
"""
|
||
if "padding" in meta:
|
||
pad = meta["padding"]
|
||
(_, end_h), (_, end_w) = pad[1:3]
|
||
h, w = image.shape[2:]
|
||
unpad_h = h - end_h
|
||
unpad_w = w - end_w
|
||
image = image[:, :, :unpad_h, :unpad_w]
|
||
image = np.clip(image / 2 + 0.5, 0, 1)
|
||
image = np.transpose(image, (0, 2, 3, 1))
|
||
# 9. Convert to PIL
|
||
if output_type == "pil":
|
||
image = self.numpy_to_pil(image)
|
||
if "src_height" in meta:
|
||
orig_height, orig_width = meta["src_height"], meta["src_width"]
|
||
image = [img.resize((orig_width, orig_height),
|
||
PIL.Image.Resampling.LANCZOS) for img in image]
|
||
else:
|
||
if "src_height" in meta:
|
||
orig_height, orig_width = meta["src_height"], meta["src_width"]
|
||
image = [cv2.resize(img, (orig_width, orig_width))
|
||
for img in image]
|
||
return image
|
||
|
||
def get_timesteps(self, num_inference_steps:int, strength:float):
|
||
"""
|
||
Helper function for getting scheduler timesteps for generation
|
||
In case of image-to-image generation, it updates number of steps according to strength
|
||
|
||
Parameters:
|
||
num_inference_steps (int):
|
||
number of inference steps for generation
|
||
strength (float):
|
||
value between 0.0 and 1.0, that controls the amount of noise that is added to the input image.
|
||
Values that approach 1.0 allow for lots of variations but will also produce images that are not semantically consistent with the input.
|
||
"""
|
||
# get the original timestep using init_timestep
|
||
init_timestep = min(int(num_inference_steps * strength), num_inference_steps)
|
||
|
||
t_start = max(num_inference_steps - init_timestep, 0)
|
||
timesteps = self.scheduler.timesteps[t_start:]
|
||
|
||
return timesteps, num_inference_steps - t_start
|
||
|
||
Configure Inference Pipeline
|
||
~~~~~~~~~~~~~~~~~~~~~~~~~~~~
|
||
|
||
First, you should create instances of OpenVINO Model.
|
||
|
||
.. code:: ipython3
|
||
|
||
from openvino.runtime import Core
|
||
core = Core()
|
||
text_enc = core.compile_model(TEXT_ENCODER_OV_PATH, "CPU")
|
||
unet_model = core.compile_model(UNET_OV_PATH, 'CPU')
|
||
vae_decoder = core.compile_model(VAE_DECODER_OV_PATH, 'CPU')
|
||
vae_encoder = core.compile_model(VAE_ENCODER_OV_PATH, 'CPU')
|
||
|
||
Model tokenizer and scheduler are also important parts of the pipeline.
|
||
Let us define them and put all components together.
|
||
|
||
.. code:: ipython3
|
||
|
||
from transformers import CLIPTokenizer
|
||
|
||
scheduler = LMSDiscreteScheduler.from_config(conf)
|
||
tokenizer = CLIPTokenizer.from_pretrained('openai/clip-vit-large-patch14')
|
||
|
||
ov_pipe = OVStableDiffusionPipeline(
|
||
tokenizer=tokenizer,
|
||
text_encoder=text_enc,
|
||
unet=unet_model,
|
||
vae_encoder=vae_encoder,
|
||
vae_decoder=vae_decoder,
|
||
scheduler=scheduler
|
||
)
|
||
|
||
Run Text-to-Image generation
|
||
~~~~~~~~~~~~~~~~~~~~~~~~~~~~
|
||
|
||
Now, you can define a text prompts for image generation and run
|
||
inference pipeline. Optionally, you can also change the random generator
|
||
seed for latent state initialization and number of steps.
|
||
|
||
.. note::
|
||
|
||
Consider increasing ``steps`` to get more precise results. A suggested value is ``50``, but it will take longer time to process.
|
||
|
||
.. code:: ipython3
|
||
|
||
import gradio as gr
|
||
from socket import gethostbyname, gethostname
|
||
|
||
|
||
def generate(prompt, negative_prompt, seed, num_steps, _=gr.Progress(track_tqdm=True)):
|
||
result = ov_pipe(
|
||
prompt,
|
||
negative_prompt=negative_prompt,
|
||
num_inference_steps=num_steps,
|
||
seed=seed,
|
||
)
|
||
return result["sample"][0]
|
||
|
||
|
||
gr.close_all()
|
||
demo = gr.Interface(
|
||
generate,
|
||
[
|
||
gr.Textbox(
|
||
"valley in the Alps at sunset, epic vista, beautiful landscape, 4k, 8k",
|
||
label="Prompt",
|
||
),
|
||
gr.Textbox(
|
||
"frames, borderline, text, charachter, duplicate, error, out of frame, watermark, low quality, ugly, deformed, blur",
|
||
label="Negative prompt",
|
||
),
|
||
gr.Slider(value=42, label="Seed", maximum=10000000),
|
||
gr.Slider(value=25, label="Steps", minimum=1, maximum=50),
|
||
],
|
||
"image",
|
||
)
|
||
ipaddr = gethostbyname(gethostname())
|
||
demo.queue().launch(server_name=ipaddr)
|